Stable Diffusion Beginner Tutorial: Generate AI Images Locally for Free
Stable Diffusion is a free, open-source AI image generator. This tutorial walks you from zero to running AI image generation on your own computer.
Why choose Stable Diffusion?
| Comparison | Stable Diffusion | Midjourney | DALL·E |
|---|---|---|---|
| Price | Free | $10–60/month | $20/month |
| Privacy | Runs locally | Cloud | Cloud |
| Control | Full control | Limited | Limited |
| Learning curve | Higher | Low | Low |
Hardware requirements
Minimum:- GPU: NVIDIA GTX 1060 6GB
- RAM: 16GB
- Disk: 20GB free space
- GPU: NVIDIA RTX 3060 12GB or better
- RAM: 32GB
- SSD: 50GB+ free space
⚠️ AMD GPUs can work too, but setup is usually more involved.
Installation
Option 1: AUTOMATIC1111 WebUI (recommended for beginners)
This is the most popular GUI.
Steps:- Install Python 3.10.6
- Install Git
- Open Terminal and run:
``
git clone https://github.com/AUTOMATIC1111/stable-diffusion-webui.git
cd stable-diffusion-webui
./webui.sh
`
- On first launch, it will download a model automatically (~4GB)
- Open http://127.0.0.1:7860 in your browser
Option 2: ComfyUI (advanced)
A node-based UI that’s more flexible—but also more complex.
Option 3: Run in the cloud
No powerful GPU? Try these services:
- Google Colab (free GPU)
- Paperspace
- RunPod
Basic usage
1. Text-to-image
Enter a description in the Prompt field.
A good prompt structure: `
subject, style, details, quality words
`
Example:
`
a beautiful sunset over mountains, oil painting style,
vibrant colors, golden hour lighting,
masterpiece, best quality, highly detailed
`
2. Key parameters
- Steps: 20–30 is a good balance
- CFG Scale: 7–12 (higher = follows prompt more strictly)
- Sampler: DPM++ 2M Karras is a solid default
- Size: 512×512 or 768×768
3. Negative prompt
Tell the model what to avoid:
`
ugly, blurry, low quality, distorted, watermark,
extra fingers, bad anatomy
`
Advanced tips
1. Use different models
Download models for different styles:
- Realistic Vision — realistic photo look
- DreamShaper — dreamy illustration
- Anything V5 — anime style
Model hub: Civitai
2. LoRA (fine-tuning add-ons)
Small add-on models to apply a specific style or character.
Usage:
Add to your prompt.
3. ControlNet
Control composition more precisely:
- Canny: edge detection
- OpenPose: human pose
- Depth: depth map
4. Inpainting
Edit part of an image:
- Upload the image
- Brush over the area you want to change
- Enter a new prompt
- Generate
Prompt writing tips
Weighting
Use brackets to adjust importance:
(important) — increase weight
((very important)) — increase more
[less important] — decrease weight
Style keywords
Photography:
`
cinematic lighting, bokeh, 35mm film,
professional photography, studio lighting
`
Painting:
`
oil painting, watercolor, digital art,
concept art, illustration
`
Anime:
`
anime style, manga, cel shading,
Studio Ghibli style
`
FAQ
Q: The hands/fingers look weird.
A: Add bad hands, extra fingers to the negative prompt, or fix it with inpainting.
Q: The image looks blurry.
A: Increase Steps to 30+, or use Hires.fix to upscale.
Q: Not enough VRAM.
A: Enable --medvram or --lowvram` in launch settings.
Recommended resources
Learning:- Stable Diffusion official docs
- Reddit r/StableDiffusion
- YouTube tutorial channels
- Civitai.com
- Hugging Face
- PromptHero
- Lexica.art
Conclusion
Stable Diffusion has a steeper learning curve, but once you get comfortable you get:
- Unlimited free generations
- Full creative freedom
- Better privacy
It’s a great investment of time.
Download AUTOMATIC1111 WebUI